Generate images with SD3.5
Image generator/identifier/reposer
Photorealistic Pron
Generate high-resolution images with text prompts
Generate images using prompts and selected LoRA models
Kolors Character to keep character developed with Flux
FLUX 4-bit Quantization(just 8GB VRAM)
High-fidelity Virtual Try-on
Generate polaroid-style images from text prompts
Easily expand image boundaries
Create customized face portraits using images and prompts
https://huggingface.co/spaces/VIDraft/mouse-webgen
Stable Diffusion 3.5 Large is an advanced AI model designed for generating high-quality images from text descriptions. It is the latest iteration in the Stable Diffusion series, building on previous versions with improved performance, better quality, and enhanced capabilities. This model is particularly noted for its ability to create detailed, realistic, and contextually relevant images based on user input.
What makes Stable Diffusion 3.5 Large different from previous versions?
Stable Diffusion 3.5 Large offers improved image quality, better context understanding, and enhanced customization options compared to earlier versions. It also includes optimizations for faster processing and more accurate results.
Do I need specialized hardware to run Stable Diffusion 3.5 Large?
Yes, running Stable Diffusion 3.5 Large requires a compatible GPU with sufficient VRAM to handle the model's computational demands. A decent NVIDIA GPU (e.g., RTX 3080 or higher) is recommended for optimal performance.
Is Stable Diffusion 3.5 Large available for public use?
Yes, Stable Diffusion 3.5 Large is available for public use, but it may require registration or access through specific platforms. Always ensure you are using the model from an official or reputable source.